How are people combining Stable Diffusion with conversational workflows?

I’ve seen more discussions lately about pairing Stable Diffusion with text-based systems, like using an AI chatbot to help refine prompts, styles, or iteration logic before image generation.
For those experimenting with this kind of setup:
Do you find conversational layers actually improve creative output, or is manual prompt tuning still better?
Interested in hearing practical experiences rather than tools or promotions

https://redd.it/1pyuowm
@rStableDiffusion

Fal has open-sourced Flux2 dev Turbo.

https://preview.redd.it/kzd4n7gs37ag1.png?width=903&format=png&auto=webp&s=ab81e3133b7d4a0922d6ebfe03296eb87b30dbeb

fal/FLUX.2-dev-Turbo · Hugging Face



https://redd.it/1pyvw0q
@rStableDiffusion

Do you think Z-Image Base release is coming soon? Recent README update looks interesting

https://redd.it/1pyszyx
@rStableDiffusion

FYI: You can train a Wan 2.2 LoRA with 16gb VRAM.

I've seen a lot of posts where people are doing initial image generation in Z-Image-Turbo and then animating it in Wan 2.2. If you're doing that solely because you prefer the aesthetics of Z-Image-Turbo, then carry on.

But for those who may be doing this out of perceived resource constraints, you may benefit from knowing that you can train LoRAs for Wan 2.2 in ostris/ai-toolkit with 16GB VRAM. Just start with the default 24GB config file and then add these parameters to your config under the model section:

layer_offloading: true
layer_offloading_text_encoder_percent: 0.6
layer_offloading_transformer_percent: 0.6

You can lower or raise the offloading percent to find what works for your setup. Of course, your batch size, gradient accumulation, and resolution all have to be reasonable as well (e.g., I did batch_size: 2, gradient_accumulation: 2, resolution: 512).

I've only tested two different LoRA runs for Wan 2.2, but so far it trains easier and, IMO, looks more natural than Z-Image-Turbo, which tends to look like it's trying to look realistic and gritty.

https://redd.it/1pz0w56
@rStableDiffusion

Qwen Image Edit 2511: Workflow for Preserving Identity & Facial Features When Using Reference Images

https://preview.redd.it/lxp8ttxre8ag1.png?width=3920&format=png&auto=webp&s=2f68e028710c494eb9a02b718696f29c8f44b4d2

Hey all,

By now many of you have experimented with the official Qwen Image Edit 2511 workflow and have run into the same issue I have: the reference image resizing inside the TextEncodeImageEditPlus node. One common workaround has been to bypass that resizing by VAE‑encoding the reference images and chaining the conditioning like:

Text Encoder → Ref Latent 1 (original) → Ref Latent 2 (ref) → Ref Latent 3 (ref)

However, when trying to transfer apparel/clothing from a reference image onto a base image, both the official workflow and the VAE‑bypass version tend to copy/paste the reference face onto the original image instead of preserving the original facial features.

I’ve been testing a different conditioning flow that has been giving me more consistent (though not perfect) results:

Text Encoder → Ref Latent 1 → Ref Latent 1 conditions Ref Latent 2 + Ref Latent 3 → combine all conditionings

From what I can tell by looking at the node code, Ref Latent 1 ends up containing conditioning from the original image and both reference images. My working theory is that re‑applying this conditioning onto the two reference latents strengthens the original image’s identity relative to the reference images.

The trade‑off is that reference identity becomes slightly weaker. For example, when transferring something like a pointed hat, the hat often “flops” instead of staying rigid—almost like gravity is being re‑applied.

I’m sure there’s a better way to preserve the base image’s identity and maintain strong reference conditioning, but I haven’t cracked it yet. I’ve also tried separately text‑encoding each image and combining them so Ref Latent 1 isn’t overloaded, but that produced some very strange outputs.

Still, I think this approach might be a step in the right direction, and maybe someone here can refine it further.

If you want to try the workflow, you can download it here:
**Pastebin Link**

Also, sampler/scheduler choice seems to matter a lot. I’ve had great results with:

er\_sde (sampler)
bong_tangent (scheduler)

(Requires the **RES4LYF** node to use these with KSampler.)

https://redd.it/1pz2gxy
@rStableDiffusion

Wan 2.2 Motion Scale - Control the Speed and Time Scale in your Wan 2.2 Videos in ComfyUI
https://youtu.be/Zmkn6_vyMN8

https://redd.it/1pz2kvv
@rStableDiffusion

[Release] I built a free, open-source desktop app to view and manage metadata (Comfy, A1111, Forge, Invoke)
https://redd.it/1pz395g
@rStableDiffusion

Flux2 Turbo Lora - Corrected ComfyUi lora keys

https://huggingface.co/ByteZSzn/Flux.2-Turbo-ComfyUI/tree/main

I converted the lora keys from **https://huggingface.co/fal/FLUX.2-dev-Turbo** to work with comfyui

https://redd.it/1pzbrg1
@rStableDiffusion

Tencent HY-Motion 1.0 - a billion-parameter text-to-motion model
https://hunyuan.tencent.com/motion?tabIndex=0

https://redd.it/1pzdjg0
@rStableDiffusion

A mysterious new year gift
https://redd.it/1pze07a
@rStableDiffusion

VNCCS V2.0 Release!

https://preview.redd.it/hl4njv7fhbag1.png?width=2979&format=png&auto=webp&s=b89fef9855370e96ebf2577ebbd5511a0e5f97b1

VNCCS - Visual Novel Character Creation Suite

VNCCS is NOT just another workflow for creating consistent characters, it is a complete pipeline for creating sprites for any purpose. It allows you to create unique characters with a consistent appearance across all images, organise them, manage emotions, clothing, poses, and conduct a full cycle of work with characters.

https://preview.redd.it/aa1gxblghbag1.png?width=2979&format=png&auto=webp&s=dd6f021a2d0da0e54030e304b98227f5393abea8

Usage

Step 1: Create a Base Character

Open the workflow VN_Step1_QWEN_CharSheetGenerator.

https://preview.redd.it/vu0yd66qhbag1.png?width=2979&format=png&auto=webp&s=8df3867624c53fec87e3a9aee29ca0d0c67b5350



# VNCCS Character Creator

First, write your character's name and click the ‘Create New Character’ button. Without this, the magic won't happen.
After that, describe your character's appearance in the appropriate fields.
SDXL is still used to generate characters. A huge number of different Loras have been released for it, and the image quality is still much higher than that of all other models.
Don't worry, if you don't want to use SDXL, you can use the following workflow. We'll get to that in a moment.

# New Poser Node

https://preview.redd.it/eigazv8shbag1.png?width=2979&format=png&auto=webp&s=6fa25b19fc6649c3b462a6bf53dc48087a773a36



# VNCCS Pose Generator

To begin with, you can use the default poses, but don't be afraid to experiment!

At the moment, the default poses are not fully optimised and may cause problems. We will fix this in future updates, and you can help us by sharing your cool presets on our [Discord](https://discord.com/invite/9Dacp4wvQw) server!

https://preview.redd.it/d52a2eduhbag1.png?width=2979&format=png&auto=webp&s=9d10d7760be342ee245863823b0089dcf72fe97c

Step 1.1 Clone any character



Try to use full body images. It can work with any images, but would "imagine" missing parst, so it can impact results.
Suit for anime and real photos

https://preview.redd.it/9bqtbpkwhbag1.png?width=3320&format=png&auto=webp&s=5f26978312ca850c60245fe7a6483f0102df7e99



# Step 2 ClothesGenerator

https://preview.redd.it/5pdwqzhyhbag1.png?width=2979&format=png&auto=webp&s=b4ace711d2c32a0641dce43296d621ab745156b7

Open the workflow `VN_Step2_QWEN_ClothesGenerator`.


Clothes helper lora are still in beta, so it can miss some "body parts" sizes. If this happens - just try again with different seeds.

# Steps 3, 4 and 5 are not changed, you can follow old guide below.

# Be creative! Now everything is possible!

https://redd.it/1pzenvt
@rStableDiffusion

Qwen Image 25-12 seen at the Horizon , Qwen Image Edit 25-11 was such a big upgrade so I am hyped
https://redd.it/1pzd675
@rStableDiffusion

VLM vs LLM prompting

https://redd.it/1pzfyq4
@rStableDiffusion

Continuous video with wan finally works!

https://reddit.com/link/1pzj0un/video/268mzny9mcag1/player

It finally happened. I dont know how a lora works this way but I'm speechless! Thanks to kijai for implementing key nodes that give us the merged latents and image outputs.
I almost gave up on wan2.2 because of multiple input was messy but here we are.

I've updated my allegedly famous workflow to implement SVI to civit AI. (I dont know why it is flagged not safe. I've always used safe examples)
https://civitai.com/models/1866565?modelVersionId=2547973

For our >!cencored!< friends;
https://pastebin.com/vk9UGJ3T

I hope you guys can enjoy it and give feedback :)

https://redd.it/1pzj0un
@rStableDiffusion